13 Best Midjourney like Stable Diffusion Checkpoints [XL and 1.5]

Researchers in artificial intelligence have long wished to create images from text instructions. This ambition is finally becoming a reality, thanks to recent developments in AI technology. Stable diffusion is an open-source AI model that can generate images from text instructions, although it has been criticized for failing to produce images of the same quality as Mid journey, a similar but non-open-source AI text-to-image model.

In this article, we will introduce you to some of the greatest stable diffusion variations that provide Mid Journey-like outcomes. Whether you’re an artist, a designer, or just someone who likes generating images, these models are sure to inspire you with their extraordinary abilities.

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OpenDalle

OpenDalle is a notable model in the realm of text-to-image AI. Developed as a checkpoint merge type, it is based on the SDXL 1.0 base model. This model stands out for its ability to adhere closely to prompts, focusing more on semantic accuracy than ultra-high-fidelity image generation. The creator, after combining several advanced models, has achieved a balance between prompt comprehension and image quality, placing OpenDalle closer to DALLE-3 in terms of understanding prompts.

The model had a minor hiccup with deletion and reuploading, and version 1.1 promises further improvements. For optimal performance, recommended settings include a CFG scale of 8 to 7, 60 to 70 steps for detailed images, and the use of DPM2 sampler and Normal or Karras scheduler.

DreamShaper XL1.0

DreamShaper XL1.0 is a powerful Stable Diffusion checkpoint that is known for its ability to generate high-quality images of people, dragons, and NSFW content. It is still under development, but it has already become a popular choice among artists and creators.

One of the biggest advantages of DreamShaper XL1.0 is that it does not require a refiner to generate high-quality images. This makes it much easier to use than other Stable Diffusion checkpoints, which can be time-consuming and difficult to set up.

DreamShaper XL1.0 is also known for its ability to generate images with sharp edges and realistic textures. This makes it ideal for creating images of people, dragons, and other complex subjects.

Additionally, DreamShaper XL1.0 is one of the best Stable Diffusion checkpoints for generating NSFW content. It can generate a wide range of NSFW images, from mild to explicit.

Crystal Clear XL

Crystal Clear XL is a powerful and versatile Stable Diffusion checkpoint that is capable of generating high-quality images in a wide range of styles. It is particularly well-suited for generating photorealistic images, but it can also be used to generate 3D, semi-realistic, and cartoonish images.

One of the key features of Crystal Clear XL is its high prompt accuracy. This means that it is able to accurately interpret your prompts and generate images that match your desired style and composition. For example, if you prompt Crystal Clear XL to generate a “close-up portrait of a young woman with blue eyes and brown hair”, it will be able to generate an image that accurately matches that description.

Crystal Clear XL also excels at generating images with realistic skin textures, eye colors, and lighting. This is due in part to the fact that it was trained on a massive dataset of high-quality images. As a result, Crystal Clear XL is able to generate images that look like they could have been taken with a professional camera.

[Lah] Mysterious | SDXL

[Lah] Mysterious | SDXL is a Stable Diffusion checkpoint that is specifically designed to generate fantasy-themed images. It is trained on a massive dataset of images and text that is related to fantasy, including historical and cyberpunk imagery, as well as data on legendary creatures.

This checkpoint is known for its ability to generate images with a unique and mysterious atmosphere. It is also capable of generating images in a variety of styles, including anime and illustration.

To generate fantasy-themed images with [Lah] Mysterious | SDXL, I recommend using a high CFG scale (3-9). This will help to produce images with a more dreamlike and atmospheric quality. You may also want to experiment with different sampling steps and sampling methods.

Illuminati Diffusion v1.1 (Unavailable)

Illuminati Diffusion v1.1 is a stable diffusion model that is similar to Midjourney in its ability to create high-quality images from text inputs. Based on stable diffusion 2.1, this model requires the use of three textural inversion embeddings to get results similar to those produced by Midjourney. The suggested safe range for generating decent results with Illuminati Diffusion v1.1 is between 768-1024 / 1152 for various seeds and prompts, and in some cases, 1280×768 or 768×1280 may also work.

Having used this model extensively, I’ve discovered that it often provides results that are more visually appealing than those generated by Midjourney. While Midjourney produces better results at times, it is fair to say that Illuminati Diffusion v1.1 is a true equal to Midjourney, not only in its ability to generate portrait images, but also in its ability to generate logo graphic designs and other image types that are typically associated with Midjourney. Overall, Illuminati Diffusion v1.1 is a robust and adaptable stable diffusion model capable of creating high-quality pictures from text inputs.

Vintedois (22h)

Vintedois (22h) is an open-source AI model designed for text-to-image generation that has garnered attention for its impressive results. The model has been trained on a large amount of high-quality images, which allows it to generate beautiful images with simple prompts without requiring a lot of prompt engineering.

One of the unique features of Vintedois (22h) is that users can enforce their preferred style by prepending their prompt with “estilovintedois.” This feature ensures that users can generate images that meet their specific style preferences. Moreover, Vintedois (22h) works well with different aspect ratios, such as 2:3 and 3:2, making it a versatile option for different use cases.

One of the strengths of Vintedois (22h) is its ability to generate high-fidelity faces with minimal steps. The model is also artistically designed, resulting in aesthetically pleasing and visually stunning images. Based on Stable Diffusion 1.5, Vintedois (22h) is also very dreamboothable, making it a popular choice for creators who want to generate high-quality images with ease.

Openjourney

Openjourney, like the Midjourney model, is a powerful and reliable Stable Diffusion model that generates high-quality images from text prompts. This model was trained using Midjourney v4 results, therefore it has learnt to make photos that roughly resemble the quality and style of Midjourney.

It is as simple to use Openjourney as it is to use any other Stable Diffusion model. You may use it through the Automatic 1111 web interface, or you can go the Prompt Hero website to get prompts that are particularly engineered to deliver the best results with Openjourney.

One of the significant benefits of Openjourney is that it is an open-source model, which means that anybody may use, modify, and fine-tune it to their own needs. As a result, it is an extremely flexible and customizable tool that can be modified to a wide range of applications.

Openjourney is also quite versatile, capable of producing images ranging from realistic landscapes and portraits to abstract art and surreal compositions. Its capacity to generate high-quality, aesthetically stunning images from basic text prompts to be an effective tool for many types of artists, designers, and content creators.

Seek Art Mega MidjrnyV4 DreamLike

Seek Art Mega MidjrnyV4 DreamLike is one of the most outstanding stable diffusion models on the market today. This model is a mixture of three separate stable diffusion models, and it can generate stunning and lifelike images from text instructions. Seek Art Mega V1, Midjrny V4, and Dreamlike Diffusion 1.0 are the three models that have been combined to form Seek Art Mega MidjrnyV4 DreamLike.

One of the most major advantages of Seek Art Mega MidjrnyV4 DreamLike is that it generates high-quality images similar to those generated by Midjourney. Seek Art Mega MidjrnyV4 DreamLike, on the other hand, is not just a copycat; it has its own distinct abilities that differentiates it apart from other versions.

This model specializes at producing vivid and vibrant images that are ideal for character and portrait illustrations. You may tweak the model to generate photos that are suited to your specific requirements. Because the model is open source, you may modify it and create your own version.

Protogen v2.2

Protogen v2.2 is a stable diffusion model that has many people amazed with its image generation capabilities. Although it was not created intentionally to resemble Midjourney, the images it generates have some of the most breathtaking compositions of any AI-generated image. In fact, it has battled with Midjourney in some situations, delivering equally outstanding results. It has even surpassed the image quality of Midjourney v4, displaying the power and versatility of Protogen v2.2.

One of Protogen v2.2’s distinctive qualities is its usage of the SD 1.5 base model, which has been trained and fine-tuned to produce visuals that are both stunning and creative. Protogen v2.2 has been praised by users for its capacity to understand the core of the prompts and transform them into graphics that are both accurate to the text and aesthetically pleasing.

Dreamlike Photoreal 2.0

Dreamlike Photoreal 2.0 is one of the most aesthetically appealing and photorealistic Stable diffusion model variants. While it is not as artistic as Midjourney, the degree of detail in its visuals is absolutely astonishing. This model is noteworthy for its versatility, since it can create images in a variety of styles, including comic style, anime style, photo-realistic, portrait, landscape, and artistic images.

One thing to keep in mind about Dreamlike Photoreal 2.0 is that it may generate NSFW images by default, thus negative prompts are needed to prevent this. Yet, the outcomes of this model are impressive, and it has the ability to produce as many artistic photos as Midjourney V4.

AyoniMix

AyoniMix is a Stable diffusion model variant that has gained popularity among AI enthusiasts and artists for its ability to produce aesthetically complex and intricate visuals. This variation has been fine-tuned to produce photos with high resolution and fine details, making it an ideal tool for artists and designers wishing to create one-of-a-kind and aesthetically appealing imagery. Although its styled output has not been widely tested, AyoniMix has been tested on several datasets and has shown its capacity to produce realistic output.

One of AyoniMix’s most notable aspects is its ability to produce a broad range of pictures, from landscapes and portraits to abstract art and surrealist imagery. This is due to the complexity of its architecture, which is meant to process multiple layers of information, such as written instructions, semantic information, and visual cues. As a result, the model may produce images that are both colorful and coherent.

Digital Diffusion – 2.1!

Digital Diffusion – 2.1 is one of the most powerful stable diffusion models for generating midjourney-like images from text instructions. This model is unique in that it is based on the stable diffusion 2.1 version, which has significant enhancements over the previous version. As a result, the model produces visuals that are fairly similar to those produced by midjourney.

One of the most distinguishing features of Digital Diffusion – 2.1 is that it is not a checkpoint model. It is a textual inversion, which means you won’t have to download enormous quantities of data to get the model. In fact, this may be downloaded for less than 100mb, making it a wonderful choice for folks who want to use stable diffusion but don’t have a lot of room to spare; nevertheless, you must first install stable diffusion 2.1.

A-Zovya RPG Artist Tools

If you are looking for a model that can create anything and everything with stunning detail and complexity, then A-Zovya RPG Artist Tools is the one for you. This general-purpose model is perfect for all your image generation needs, no matter how intricate or challenging they may be.

A-Zovya RPG Artist Tools comes with several models that cater to your specific requirements. For instance, if you want to generate pretty girl portraits, you can use the SD 1.5 models for the best results. These models come with a stronger painterly style, higher contrast, and sharper images that showcase more RPG knowledge. If you need better coherency on non-human objects or hands with only five fingers, then the SD 2.1 models are perfect for you. These models offer more details and a higher native resolution, making them ideal for creating complex and intricate images.

A-Zovya RPG Artist Tools offers a default aesthetic for each of its models, making it easier for users to select the style they want. You can also prompt any style you need with these models, making them incredibly versatile and flexible. Additionally, the model comes with a detailed tutorial on how to get the results you see in the preview images, so you can achieve the same stunning images with ease.

The Stable Diffusion 1.5 (512) versions of A-Zovya RPG Artist Tools come with several different options to choose from. The V2 Stronger painterly style offers a higher contrast and sharper images, while the V2 offset Noise Offset adds more contrast and brings the model back to photorealism. The V2 Art Trained model is perfect for those who love the digital painting style pre-AI, as it comes with fewer details and bigger brush strokes to mimic that style. The V2 inpaint model is excellent for outpainting, while the V1 Smoother renders offer the least painterly effect. Finally, the V1 inpaint model is perfect for outpainting as well.

If you’re looking for even more detail and resolution, the Stable Diffusion 2.1 (768) versions of A-Zovya RPG Artist Tools are perfect for you. The SD 2.1 768 V1 model offers a strong painterly style that is very coherent with hands and objects. However, it’s not suitable for nudity, so keep that in mind when selecting this model.

Perpetual Diffusion 1.0

Perpetual Diffusion 1.0 is a versatile and powerful generative model that uses stable diffusion to generate high-quality images. Based on the stable diffusion SD 2.1 768 model, Perpetual Diffusion 1.0 is specifically trained to output 1024×1024 images in a square format. It also accepts large ratios (panoramas) with less duplication than the vanilla SD 2.1 model.

One of the unique features of Perpetual Diffusion 1.0 is its ability to create a wide variety of images, thanks to its general-purpose nature. Whether you want to generate landscapes, portraits, or abstract art, Perpetual Diffusion 1.0 can do it all. Additionally, there are two versions of the model available: Perpetual Sun, which creates brighter images, and Perpetual Moon, which creates darker images.

To get the best results from Perpetual Diffusion 1.0, it is recommended to use certain types of negative and positive prompts. Negative prompts such as “bad art, poorly drawn, blur, low quality, smooth, mutated, jpg artifacts, horror, overexposed, oversaturated, undersaturated, amateur” can help the model avoid these undesirable characteristics in its generated images. Positive prompts such as “4k, high quality, very detailed” can help the model produce more realistic and detailed images.

Perpetual Diffusion 1.0 is also compatible with the ControlNet for SD 2.1, allowing users to have even more control over the generated images. The model can be hosted on non-commercial websites or digital platforms, but commercial use requires prior written consent from the authors.